Text Guided Real Image Editing Via Diffusion Model Inversion

via “fast image generation with distilled diffusion steps”

Stability AI's 8B parameter flagship image generation model.

Unique: Applies knowledge distillation to compress diffusion steps from standard schedule to 4 steps while preserving the full 8.1B parameter model, enabling faster inference without architectural changes or separate lightweight model training

vs others: Faster than standard Stable Diffusion 3.5 Large with same parameter count, but slower than purpose-built fast models like LCM-LoRA or consistency models; trades speed for quality more conservatively than extreme distillation approaches

via “text-to-image generation with diffusion model inference”

Invoke is a leading creative engine for Stable Diffusion models, empowering professionals, artists, and enthusiasts to generate and create visual media using the latest AI-driven technologies. The solution offers an industry leading WebUI, and serves as the foundation for multiple commercial product

Unique: Uses a node-based invocation graph architecture (BaseInvocation system) that decouples model inference from UI, enabling reusable, composable generation pipelines where each step (conditioning, sampling, post-processing) is a discrete node with schema-driven validation and serialization. This contrasts with monolithic pipeline approaches by allowing users to visually construct custom workflows.

vs others: Offers more granular control over generation parameters and pipeline composition than consumer tools like Midjourney, while maintaining ease-of-use through a professional WebUI; faster iteration than cloud APIs due to local model execution and no network latency.

via “single-step text-to-image generation with latency optimization”

text-to-image model by undefined. 13,26,546 downloads.

Unique: Implements single-step diffusion via knowledge distillation from larger teacher models, collapsing 20-50 sampling iterations into one forward pass while maintaining competitive image quality — a fundamentally different architecture from iterative refinement models like SDXL that require sequential denoising steps

vs others: Achieves 10-50x faster inference than SDXL or Flux with comparable quality on standard prompts, making it the fastest open-source text-to-image model for latency-critical applications, though with trade-offs in detail complexity and style control

via “single-step text-to-image generation with adversarial diffusion distillation”

text-to-image model by undefined. 8,95,582 downloads.

Unique: Uses adversarial diffusion distillation (ADD) to compress SDXL's 50-step inference into a single forward pass, achieving ~40× speedup while maintaining competitive image quality through adversarial training against a discriminator that enforces perceptual similarity to multi-step outputs.

vs others: 40× faster than standard SDXL 1.0 (0.5s vs 20s on RTX 3090) while maintaining comparable aesthetic quality, making it the only open-source text-to-image model suitable for real-time interactive applications without sacrificing photorealism.

via “text-to-image generation”

text-to-image model by undefined. 2,75,100 downloads.

Unique: Utilizes a refined latent diffusion approach that balances quality and computational efficiency, allowing for faster image generation compared to earlier iterations.

vs others: Generates images with higher fidelity and detail than previous models like Stable Diffusion 2.1, thanks to improved training techniques and dataset diversity.

via “diffusion-based iterative image synthesis with guidance”

text-to-image model by undefined. 3,26,804 downloads.

Unique: Implements diffusion-based synthesis as a core capability rather than relying on external diffusion frameworks, with integrated guidance mechanism that balances prompt adherence against image quality through learned weighting of conditional and unconditional predictions

vs others: More flexible than GAN-based approaches (single-step generation) by enabling mid-generation adjustments through guidance, and more efficient than autoregressive pixel-space models by operating in compressed latent space

via “image inpainting”

Stable Diffusion by Stability AI is a state of the art text-to-image model that generates images from text. #opensource

Unique: The inpainting feature is integrated into the same diffusion process as the text-to-image generation, allowing for a unified model that can handle both tasks without needing separate architectures.

vs others: More flexible than traditional inpainting tools because it can generate entirely new content based on textual prompts rather than relying solely on existing image data.

via “inpainting-selective-image-region-replacement”

Diffusion Bee is the easiest way to run Stable Diffusion locally on your M1 Mac. Comes with a one-click installer. No dependencies or technical knowledge needed.

Unique: Uses specialized inpainting model checkpoints that are trained with mask-aware conditioning, allowing the diffusion process to understand mask boundaries and blend seamlessly. The implementation encodes both image and mask through separate pathways in the latent space, enabling precise control over which regions are modified.

vs others: More precise than content-aware fill algorithms (which use statistical inpainting) and faster than manual Photoshop cloning, while requiring less training data than generative inpainting models that must learn from scratch.

via “decomposed dual-branch diffusion inpainting with masked feature separation”

[ECCV 2024] The official implementation of paper "BrushNet: A Plug-and-Play Image Inpainting Model with Decomposed Dual-Branch Diffusion"

Unique: Uses decomposed dual-branch architecture with dense per-pixel control injected at multiple UNet resolution levels, enabling plug-and-play integration without modifying base model weights. Unlike naive masking approaches, separates masked feature processing from latent noise processing, reducing learning burden and improving boundary quality.

vs others: Achieves higher inpainting quality than simple mask-based approaches (e.g., Inpaint-LoRA) while maintaining compatibility with any pre-trained diffusion model, and requires significantly less training data than full model fine-tuning approaches.

via “image-to-image transformation with text-guided refinement”

Kandinsky 2 — multilingual text2image latent diffusion model

Unique: Uses MOVQ encoder (67M parameters) instead of standard VAE for input image encoding, providing better reconstruction fidelity in latent space. Strength parameter controls noise schedule initialization, enabling smooth interpolation between preservation and regeneration without separate model variants.

vs others: Achieves finer control over image preservation than Stable Diffusion's img2img through explicit diffusion prior conditioning, and supports multilingual prompts natively unlike most open-source alternatives.

via “differential diffusion with region-specific generation control”

我的 ComfyUI 工作流合集 | My ComfyUI workflows collection

Unique: Provides differential diffusion workflows that expose per-pixel generation strength control, a capability unavailable in most commercial tools (Midjourney, DALL-E 3) and rarely documented in open-source implementations

vs others: More granular than inpainting masks (binary or soft) because differential diffusion allows continuous per-pixel strength variation; more flexible than ControlNet because it operates on the image itself rather than requiring separate control images

via “practical stable diffusion applications (inpainting, editing, upscaling)”

Python materials for the online course on diffusion models by [@huggingface](https://github.com/huggingface).

via “instruction-guided image editing via diffusion”

instruct-pix2pix — AI demo on HuggingFace

Unique: Uses a dual-conditioning architecture combining CLIP text embeddings with image features in a single UNet, enabling instruction-guided edits without separate mask inputs or region selection — differs from traditional inpainting approaches that require explicit mask specification

vs others: More intuitive than mask-based editing tools and faster than training custom LoRA adapters, but less precise than pixel-level editing tools like Photoshop for geometric transformations

via “text-to-image generation with diffusion-based synthesis”

IF — AI demo on HuggingFace

Unique: Implements a cascaded multi-stage diffusion pipeline (base + super-resolution stages) rather than single-stage generation, enabling higher quality and resolution through progressive refinement. Uses frozen language model embeddings for text conditioning, reducing training complexity compared to end-to-end approaches like DALL-E.

vs others: Achieves higher image quality and finer detail than single-stage models (Stable Diffusion) through cascaded architecture, while maintaining faster inference than autoregressive approaches (DALL-E) by leveraging efficient diffusion sampling.

via “text-to-image generation within masked regions using diffusion models”

MagicQuill — AI demo on HuggingFace

Unique: Integrates text-conditioned diffusion inpainting via a pre-trained model hosted on HuggingFace, eliminating the need for local GPU setup. The Gradio interface abstracts model loading, tokenization, and inference orchestration into a simple prompt-and-mask input flow.

vs others: More accessible than running Stable Diffusion locally because it requires no GPU or software installation, though with less control over advanced parameters (guidance scale, scheduler, negative prompts) than command-line tools like Automatic1111.

via “diffusion model inference with gpu acceleration”

IC-Light — AI demo on HuggingFace

Unique: Implements lighting-aware conditioning by injecting spatial maps into the diffusion model's cross-attention layers, rather than relying solely on text prompts or implicit context. This allows precise control over lighting direction without requiring complex prompt engineering.

vs others: Faster than CPU-based inference by 50-100x due to GPU parallelization of matrix operations, and produces higher-quality results than simpler inpainting methods (like content-aware fill) because it leverages learned generative priors from large-scale training.

via “prompt-to-image inference with model selection”

Z-Image-Turbo — AI demo on HuggingFace

Unique: Model selection is implemented as Gradio UI components bound directly to HuggingFace Inference API model identifiers, allowing runtime model switching without backend code changes — the Space configuration itself defines available models

vs others: Simpler than ComfyUI for model comparison because it abstracts away node graphs and requires no local VRAM, but less flexible than Ollama for fine-grained model parameter control

via “text-guided image editing with minimal denoising steps”

* ⭐ 10/2022: [LAION-5B: An open large-scale dataset for training next generation image-text models (LAION-5B)](https://arxiv.org/abs/2210.08402)

Unique: Achieves 2-4 step image editing by distilling guidance information, enabling interactive editing without separate guidance models. Preserves unedited regions through latent-space conditioning while reducing computational overhead.

vs others: 10-50× faster than standard diffusion-based editing (e.g., InstructPix2Pix with full steps), but may sacrifice fine-grained control and semantic accuracy compared to non-distilled approaches.

via “text-to-image generation with diffusion model inference”

IllusionDiffusion — AI demo on HuggingFace

Unique: Integrates optical illusion conditioning into the standard Stable Diffusion pipeline via cross-attention fusion, rather than using simple prompt engineering or post-processing, enabling structural guidance that persists throughout the entire denoising process

vs others: Produces more coherent illusion-guided outputs than naive prompt-based approaches because the illusion pattern is embedded directly into the diffusion latent space, not just mentioned in text; faster than fine-tuning custom models because it uses pre-trained Stable Diffusion weights with conditioning injection

via “image-inpainting-via-conditional-diffusion”

* 🏆 2020: [An Image is Worth 16x16 Words: Transformers for Image Recognition at Scale (ViT)](https://arxiv.org/abs/2010.11929)

Unique: DDPM enables zero-shot inpainting by leveraging the forward process to compute noisy versions of known pixels at each timestep, then replacing unknown pixels with model predictions. This approach requires no special training and works with any trained diffusion model. The key insight is that the forward process provides a principled way to inject known information at each denoising step.

vs others: Requires no special training (unlike GAN-based inpainting), enables flexible mask shapes and sizes, and can be combined with text guidance for semantic inpainting.